Stable Diffusion Illusion

Illusion Diffusion converts uploaded photos into optical‑illusion artwork using Stable Diffusion and ControlNet. Users set prompts and Illusion Strength, adjust advanced controls, and select patterns to produce high‑quality, shareable visual transformations quickly.

Stable Diffusion Online lets users generate photo‑realistic images from text using the Stable Diffusion XL model. It offers fast GPU‑accelerated rendering, real‑time inpainting/outpainting, a 9‑million‑entry prompt database, and no prompt or image storage.

Think Diffusion is an AI art tool that brings pro-level AI art capabilities to anyone with a browser in just a few clicks, providing top-tier models and ControlNet, and allowing for easy launch of additional virtual machines for simultaneous use.

Stable Diffusion 3 is a text‑to‑image model using Diffusion Transformer architecture, offering accurate multi‑subject rendering, high‑detail outputs, noise‑prediction sampling, safety filters, API and downloadable weights for scalable, local or integrated deployment.

Stabl Diffus is an AI tool for generating images from text prompts using a stable image generation model suitable for government and creative ML usage.

Unstable Diffusion is an AI image generator that provides unrestricted and uncensored image generation from text prompts.

The Stable Diffusion AI Generator is an online platform with a 6 billion parameter dual model system and simplified prompting, allowing users to easily generate stunning visuals from text using Imagine AI art generator.

Provides API access to pretrained image generation models for text‑to‑image, image‑to‑image, and inpainting, with real‑time editing. Supports single‑call Dreambooth/LoRA training without local GPU, plus voice cloning, text‑to‑3D, interior design, and video creation.

DiffusionBee runs locally on macOS and Intel‑based Macs, offering Stable Diffusion text‑to‑image, upscaling, variant creation, image‑to‑image, generative fill, video tools, an AI canvas for iterative editing, and custom model training—All on‑device for privacy.

A free, user-friendly, multilingual, and open-source AI image generator that utilizes Stable Diffusion.

AI art generator in cloud

This AI tool generates images and text using machine learning features and has built-in safety measures, with ongoing development and a Discord server for help and support.

Happy Diffusion runs Stable Diffusion in the browser, enabling instant adult image creation with 50+ pre‑integrated models and unlimited Civitai models. It uses an NVIDIA A100 GPU, handles up to 7,000 images/hour, and erases data per session.

QR Diffusion blends Stable Diffusion and ControlNet to create customizable QR codes with artistic imagery. Users can edit codes, embed long URLs, track scan analytics, and overlay designs via an insert feature for brand‑consistent graphics.

Stability AI has recently launched the Stable Animation SDK, a text-to-animation tool for developers that enables advanced stable diffusion models.

Finetuned Diffusion is a pre-trained model for text generation tasks on Hugging Face's Space platform. It can be used in Colab without waiting in a queue.

AnimateDiff transforms text prompts or static images into short video clips by generating key frames with Stable Diffusion and interpolating realistic motion. It supports ControlNet editing, works in AUTOMATIC1111 Web UI and Colab, and requires an 8 GB GPU.

Scribble Diffusion AI transforms hand‑drawn sketches into polished images with one click, allowing browser drawing or file upload. It offers prompt‑based style variations, quick refinement, and export in common formats for seamless pipeline integration.

Story Diffusion is an AI tool for generating comic stories with consistent character portraits. Users can upload photos, choose from six storytelling styles, and create engaging comic content while ensuring data privacy through encryption.

StoryDiffusion generates consistent comic panels from text prompts, preserving character likeness via reference uploads. It offers high‑quality layouts, flexible styles, video conversion with narration, and secure export options for writers, illustrators, and educators.

Image Pipeline delivers AI image creation and editing using Stable Diffusion, Flux, and custom checkpoints. It supports LoRA, embeddings, adapters, ControlNet for inpainting, and Face Lock/Quick Swap for facial editing, all via a REST API.

MagicAnimate animates static images with temporally consistent motion using a diffusion model and MSE‑finetuned VAE, supports cross‑ID animation, pose‑driven motion via OpenPose, and can animate text‑generated images. Open‑source, local or via Hugging Face/Replicate.

Flying Dog plugs Stable Diffusion, DALL‑E 2, and Automatic1111 into Photoshop, offering Text‑to‑Image, Image‑to‑Image, Inpainting, and Outpainting on layers. It runs on GPU‑accelerated local servers (≥6 GB NVIDIA) or cloud for macOS, replacing stock searches and Content‑Aware Fill with precise visua

AI Disturbance Overlay uses adversarial perturbations to add invisible pixel changes that disrupt AI image recognition and style extraction. It protects artwork against generative models, preserving artists’ creative ownership in the browser.

Web-based Dreambooth Google Colab and web UI interface for stable Diffusion.

Story Diffusion Gen is a digital storytelling platform that generates high-quality images and videos from text prompts. It offers tools for comic layout, visual continuity, and motion prediction, catering to creators of all skill levels.

Stable Video Diffusion Online converts still images and text prompts into 4‑second, watermark‑free video clips in 16:9, 9:16, or 1:1 aspect ratios. Fast, server‑based processing supports media, advertising, and education workflows, with open‑source code for customization.

Flux AI Image Generator creates high‑quality images in about 30 seconds using just 1–4 diffusion steps. It supports 0.1–2 MP resolutions, varied aspect ratios, styles, and artistic modes. Three variants (Pro, Schnell, Dev) balance speed and quality.

SV3D transforms a single image into multi‑angled views and 3D meshes using Stable Video Diffusion. It offers free‑to‑download weights for researchers, supports orbital or camera‑path videos, and aids rapid 3D asset creation for designers and developers.

SDXL Turbo is a text‑to‑image model using Adversarial Diffusion Distillation for single‑step, high‑quality 512×512 outputs in under a second on modern GPUs. It supports multiple text encoders, is open‑source, and fits real‑time applications.

DiffusionArt converts text prompts into high-resolution images using diffusion models, offering prompt-to-image, inpainting, image editing, batch generation, upscaling and style controls for consistent, print- and web-ready visuals with workflow and access management.

ComfyUI is a powerful and modular stable diffusion GUI with a graph/nodes interface that lets you design and execute advanced stable diffusion pipelines using a flowchart-based interface.

Story Diffusion AI generates narrative content from prompts in styles like classic literature, adventure, or contemporary prose. Users can adjust tone, setting, and character voice with controls, and the platform produces coherent visual scenes using self‑attention continuity across long sequences.

The Stable Diffusion AI tool provides 20 AI models for download on Mac through their unofficial Telegram bot and offers support via the @stable_diffusion_bot or @doeventchat channel. The tool can also send messages via Telegram using the AI models.

Fashion Diffusion AI is an all-in-one generative AI platform for fashion design, providing tools for concept generation, virtual try-on, and asset creation. It streamlines the entire design-to-market workflow, enabling rapid prototyping, fit evaluation, and production of marketing visuals.